🎨 Image generation

(Generated with AnimagineXL)

LocalAI supports generating images with Stable diffusion, running on CPU using C++ and Python implementations.

Usage

OpenAI docs: https://platform.openai.com/docs/api-reference/images/create

To generate an image you can send a POST request to the /v1/images/generations endpoint with the instruction as the request body:

curl http://localhost:8080/v1/images/generations -H "Content-Type: application/json" -d '{
  "prompt": "A cute baby sea otter",
  "size": "256x256"
}'

Available additional parameters: mode, step.

Note: To set a negative prompt, you can split the prompt with |, for instance: a cute baby sea otter|malformed.

curl http://localhost:8080/v1/images/generations -H "Content-Type: application/json" -d '{
  "prompt": "floating hair, portrait, ((loli)), ((one girl)), cute face, hidden hands, asymmetrical bangs, beautiful detailed eyes, eye shadow, hair ornament, ribbons, bowties, buttons, pleated skirt, (((masterpiece))), ((best quality)), colorful|((part of the head)), ((((mutated hands and fingers)))), deformed, blurry, bad anatomy, disfigured, poorly drawn face, mutation, mutated, extra limb, ugly, poorly drawn hands, missing limb, blurry, floating limbs, disconnected limbs, malformed hands, blur, out of focus, long neck, long body, Octane renderer, lowres, bad anatomy, bad hands, text",
  "size": "256x256"
}'

Backends

stablediffusion-ggml

This backend is based on stable-diffusion.cpp. Every model supported by that backend is supported indeed with LocalAI.

Setup

There are already several models in the gallery that are available to install and get up and running with this backend, you can for example run flux by searching it in the Model gallery (flux.1-dev-ggml) or start LocalAI with run:

local-ai run flux.1-dev-ggml

To use a custom model, you can follow these steps:

Create a model file stablediffusion.yaml in the models folder:

name: stablediffusion
backend: stablediffusion-ggml
parameters:
  model: gguf_model.gguf
step: 25
cfg_scale: 4.5
options:
- "clip_l_path:clip_l.safetensors"
- "clip_g_path:clip_g.safetensors"
- "t5xxl_path:t5xxl-Q5_0.gguf"
- "sampler:euler"

Download the required assets to the models repository
Start LocalAI

Diffusers

Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. LocalAI has a diffusers backend which allows image generation using the diffusers library.

(Generated with AnimagineXL)

Model setup

The models will be downloaded the first time you use the backend from huggingface automatically.

Create a model configuration file in the models directory, for instance to use Linaqruf/animagine-xl with CPU:

name: animagine-xl
parameters:
  model: Linaqruf/animagine-xl
backend: diffusers

f16: false
diffusers:
  cuda: false # Enable for GPU usage (CUDA)
  scheduler_type: euler_a

Dependencies

This is an extra backend - in the container is already available and there is nothing to do for the setup. Do not use core images (ending with -core). If you are building manually, see the build instructions.

Model setup

The models will be downloaded the first time you use the backend from huggingface automatically.

Create a model configuration file in the models directory, for instance to use Linaqruf/animagine-xl with CPU:

name: animagine-xl
parameters:
  model: Linaqruf/animagine-xl
backend: diffusers
cuda: true
f16: true
diffusers:
  scheduler_type: euler_a

Local models

You can also use local models, or modify some parameters like clip_skip, scheduler_type, for instance:

name: stablediffusion
parameters:
  model: toonyou_beta6.safetensors
backend: diffusers
step: 30
f16: true
cuda: true
diffusers:
  pipeline_type: StableDiffusionPipeline
  enable_parameters: "negative_prompt,num_inference_steps,clip_skip"
  scheduler_type: "k_dpmpp_sde"
  clip_skip: 11

cfg_scale: 8

Configuration parameters

The following parameters are available in the configuration file:

Parameter	Description	Default
`f16`	Force the usage of `float16` instead of `float32`	`false`
`step`	Number of steps to run the model for	`30`
`cuda`	Enable CUDA acceleration	`false`
`enable_parameters`	Parameters to enable for the model	`negative_prompt,num_inference_steps,clip_skip`
`scheduler_type`	Scheduler type	`k_dpp_sde`
`cfg_scale`	Configuration scale	`8`
`clip_skip`	Clip skip	None
`pipeline_type`	Pipeline type	`AutoPipelineForText2Image`
`lora_adapters`	A list of lora adapters (file names relative to model directory) to apply	None
`lora_scales`	A list of lora scales (floats) to apply	None

There are available several types of schedulers:

Scheduler	Description
`ddim`	DDIM
`pndm`	PNDM
`heun`	Heun
`unipc`	UniPC
`euler`	Euler
`euler_a`	Euler a
`lms`	LMS
`k_lms`	LMS Karras
`dpm_2`	DPM2
`k_dpm_2`	DPM2 Karras
`dpm_2_a`	DPM2 a
`k_dpm_2_a`	DPM2 a Karras
`dpmpp_2m`	DPM++ 2M
`k_dpmpp_2m`	DPM++ 2M Karras
`dpmpp_sde`	DPM++ SDE
`k_dpmpp_sde`	DPM++ SDE Karras
`dpmpp_2m_sde`	DPM++ 2M SDE
`k_dpmpp_2m_sde`	DPM++ 2M SDE Karras

Pipelines types available:

Pipeline type	Description
`StableDiffusionPipeline`	Stable diffusion pipeline
`StableDiffusionImg2ImgPipeline`	Stable diffusion image to image pipeline
`StableDiffusionDepth2ImgPipeline`	Stable diffusion depth to image pipeline
`DiffusionPipeline`	Diffusion pipeline
`StableDiffusionXLPipeline`	Stable diffusion XL pipeline
`StableVideoDiffusionPipeline`	Stable video diffusion pipeline
`AutoPipelineForText2Image`	Automatic detection pipeline for text to image
`VideoDiffusionPipeline`	Video diffusion pipeline
`StableDiffusion3Pipeline`	Stable diffusion 3 pipeline
`FluxPipeline`	Flux pipeline
`FluxTransformer2DModel`	Flux transformer 2D model
`SanaPipeline`	Sana pipeline

Advanced: Additional parameters

Additional arbitrarly parameters can be specified in the option field in key/value separated by ::

name: animagine-xl
options:
- "cfg_scale:6"

Note: There is no complete parameter list. Any parameter can be passed arbitrarly and is passed to the model directly as argument to the pipeline. Different pipelines/implementations support different parameters.

The example above, will result in the following python code when generating images:

pipe(
    prompt="A cute baby sea otter", # Options passed via API
    size="256x256", # Options passed via API
    cfg_scale=6 # Additional parameter passed via configuration file
)

Usage

Text to Image

Use the image generation endpoint with the model name from the configuration file:

curl http://localhost:8080/v1/images/generations \
    -H "Content-Type: application/json" \
    -d '{
      "prompt": "<positive prompt>|<negative prompt>", 
      "model": "animagine-xl", 
      "step": 51,
      "size": "1024x1024" 
    }'

Image to Image

https://huggingface.co/docs/diffusers/using-diffusers/img2img

An example model (GPU):

name: stablediffusion-edit
parameters:
  model: nitrosocke/Ghibli-Diffusion
backend: diffusers
step: 25
cuda: true
f16: true
diffusers:
  pipeline_type: StableDiffusionImg2ImgPipeline
  enable_parameters: "negative_prompt,num_inference_steps,image"

IMAGE_PATH=/path/to/your/image
(echo -n '{"file": "'; base64 $IMAGE_PATH; echo '", "prompt": "a sky background","size": "512x512","model":"stablediffusion-edit"}') |
curl -H "Content-Type: application/json" -d @-  http://localhost:8080/v1/images/generations

🖼️ Flux kontext with `stable-diffusion.cpp`

LocalAI supports Flux Kontext and can be used to edit images via the API:

Install with:

local-ai run flux.1-kontext-dev

To test:

curl http://localhost:8080/v1/images/generations -H "Content-Type: application/json" -d '{
  "model": "flux.1-kontext-dev",
  "prompt": "change 'flux.cpp' to 'LocalAI'",
  "size": "256x256",
  "ref_images": [
  	"https://raw.githubusercontent.com/leejet/stable-diffusion.cpp/master/assets/flux/flux1-dev-q8_0.png"
  ]
}'

Depth to Image

https://huggingface.co/docs/diffusers/using-diffusers/depth2img

name: stablediffusion-depth
parameters:
  model: stabilityai/stable-diffusion-2-depth
backend: diffusers
step: 50
f16: true
cuda: true
diffusers:
  pipeline_type: StableDiffusionDepth2ImgPipeline
  enable_parameters: "negative_prompt,num_inference_steps,image"

cfg_scale: 6

(echo -n '{"file": "'; base64 ~/path/to/image.jpeg; echo '", "prompt": "a sky background","size": "512x512","model":"stablediffusion-depth"}') |
curl -H "Content-Type: application/json" -d @-  http://localhost:8080/v1/images/generations

img2vid

name: img2vid
parameters:
  model: stabilityai/stable-video-diffusion-img2vid
backend: diffusers
step: 25
f16: true
cuda: true
diffusers:
  pipeline_type: StableVideoDiffusionPipeline

(echo -n '{"file": "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/svd/rocket.png?download=true","size": "512x512","model":"img2vid"}') |
curl -H "Content-Type: application/json" -X POST -d @- http://localhost:8080/v1/images/generations

txt2vid

name: txt2vid
parameters:
  model: damo-vilab/text-to-video-ms-1.7b
backend: diffusers
step: 25
f16: true
cuda: true
diffusers:
  pipeline_type: VideoDiffusionPipeline
  cuda: true

(echo -n '{"prompt": "spiderman surfing","size": "512x512","model":"txt2vid"}') |
curl -H "Content-Type: application/json" -X POST -d @- http://localhost:8080/v1/images/generations

🎨 Image generation

Usage

Backends

stablediffusion-ggml

Setup

Diffusers

Model setup

Dependencies

Model setup

Local models

Configuration parameters

Advanced: Additional parameters

Usage

Text to Image

Image to Image

🖼️ Flux kontext with stable-diffusion.cpp

Depth to Image

img2vid

txt2vid

🖼️ Flux kontext with `stable-diffusion.cpp`